Jack000 / glid-3-xl-stable
stable diffusion training
☆291Updated 2 years ago
Related projects ⓘ
Alternatives and complementary repositories for glid-3-xl-stable
- Discord bot and Interface for Stable Diffusion☆281Updated last year
- ☆242Updated 2 years ago
- Implementation of Paint-with-words with Stable Diffusion : method from eDiff-I that let you generate image from text-labeled segmentation…☆636Updated last year
- Advanced fine tuning tools for vision models☆226Updated last year
- 1.4B latent diffusion model fine tuning☆261Updated 2 years ago
- Fast finetuning using a booster model that puts the initial state to a local minimum☆113Updated last year
- Implementation of Encoder-based Domain Tuning for Fast Personalization of Text-to-Image Models☆321Updated last year
- Majesty Diffusion by @Dango233(@Dango233max) and @apolinario (@multimodalart)☆276Updated 2 years ago
- [ECCV 2022] Compositional Generation using Diffusion Models☆455Updated 3 months ago
- This project helps you do prompt-based inpainting without having to paint the mask - using Stable Diffusion and Clipseg☆368Updated 2 years ago
- General fine tuning for Stable Diffusion☆506Updated last year
- Mixture of Diffusers for scene composition and high resolution image generation☆418Updated last year
- Dreambooth implementation based on Stable Diffusion with minimal code.☆142Updated 2 years ago
- Adaptation of the merging method described in the paper - Git Re-Basin: Merging Models modulo Permutation Symmetries (https://arxiv.org/a…☆139Updated 6 months ago
- ☆165Updated 2 years ago
- ☆110Updated 2 years ago
- Home of `erlich` and `ongo`. Finetune latent-diffusion/glid-3-xl text2image on your own data.☆182Updated 2 years ago
- ☆318Updated 2 months ago
- adaptation of huggingface's dreambooth training script to support depth2img☆101Updated last year
- ELITE: Encoding Visual Concepts into Textual Embeddings for Customized Text-to-Image Generation (ICCV 2023, Oral)☆514Updated 10 months ago
- Frontend for deeplearning Image generation☆144Updated 5 months ago
- Create butter-smooth transitions between prompts, powered by stable diffusion☆355Updated 7 months ago
- Dataset of prompts, synthetic AI generated images, and aesthetic ratings.☆399Updated 2 years ago
- Code for Enhancing Detail Preservation for Customized Text-to-Image Generation: A Regularization-Free Approach☆466Updated 10 months ago
- An implementation of a server for the Stability AI Stable Diffusion API☆174Updated last year
- Implementation of Dreambooth (https://arxiv.org/abs/2208.12242) with Stable Diffusion (tweaks focused on training faces)☆143Updated last year
- ☆157Updated last year
- Improve the editability of any Stability Diffusion subject while retaining a high degree of likeness☆150Updated 8 months ago
- ☆204Updated 2 years ago
- Finetuning SD in style.☆670Updated last year